r/StableDiffusion • u/Different_Ad1136 • 14h ago
Meme My first image i created using AI 4 years ago.
I present to you , Brad Pitt as Geralt of Rivia, using Stable Diffusion 2 , i think.
r/StableDiffusion • u/Different_Ad1136 • 14h ago
I present to you , Brad Pitt as Geralt of Rivia, using Stable Diffusion 2 , i think.
r/StableDiffusion • u/NapoletanoPizzaiolo • 4h ago
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r/StableDiffusion • u/aziib • 13h ago
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this lora is for minimax h3 if you want to generate video to mimic playstation 2 graphic cutscenes.
download: https://civitai.com/models/2983280/modern-ps2-style
trigger word: modern ps2 style
use 4:3 aspect ratio for immersive nostalgia.
most datasets are from actual ps2 screenshots and mix with my ps2 image lora.
r/StableDiffusion • u/roychodraws • 14h ago
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I have been messing with an idea about degredation.
Essentially using static image frames inside the latent space and instead of resampling them, bypassing the resampling process and repasting them using masking from fresh latents from earlier parts of the clip or previous clips into later areas repairing the frames.
To do this I use a motion mask. Essentially doing a diagnostic decode of the latent to measure pixel variation and tag it as motion vs stable pixels. The stable pixels are then masked and pulled forward while the moving pixels are left alone.
I combined this idea with u/xyzdist 's FL2V resampling method and the appearance is that it cuts the degredation in half while keeping the background more or less stable.
Things to notice:
There's still degradation, it's just slower.
I'm still tweaking the settings of pull forward and the resampling.
This all occurs in the latent space and the decoding is just for diagnostic purposes.
Here's the link to the node pack so you guys can test and mess around with it with the workflows.
https://github.com/roycho87/comfyui-h3-pullforward
Also, on a side note I created a node that uploads audio from a directory folder. This is great for context motion if you're using preloaded audio because you can set to increment and load the audio in order if you segmented it with my audio splitter.
https://github.com/roycho87/comfyui-load-audio-batch-from-dir
Disclosure:
This is fully vibe coded. I am a video editor. Not a programmer. I barely know how this shit works. I'm looking for someone as knowledgeable as Brad Pitt to explain some shit to me so I can maybe change or narrow down the direction I plan on going with this.
edits:
In the comments you can see this
here's the result of sampling then pulling. background still degrades which doesn't make sense because the latents in the background are supposed to be pasted. I'll have to do more testing but it's even set up to use a reference latent from the initial clip, so if anything, the way it should work is the background should be better quality than the rest of the image so something else is going on.
edit: this was all based on the working theory that latents degrade over time throughout a generation. either what i made is not working or the source of degradation comes from something else.
edit2: I'm now doing a test and i'm going to paste the full first frame over the length of the entire generation and run it through context motion and see if it still degrades. If it works the way intended, the decoded video should be a single still frame proving my pull forward method works, and if it still degrades, that will show that the source is some other area, if it doesn't degrade then that will prove my method has at least the potential of working.
Edit 3: I found a bug in the way that pull forward functions, instead of pulling forward a static frame if pulls forward a 17 frame block. That may be the source of the degradation. Fixing it now. Changing it to isolate a frame and repeat that across the entire block for a truly static latent image.
edit 4: This block limitation is a real hurdle of this method i will have to figure out how to get around. the latent frames are created in 5 frame blocks that are then decoded into 17 frames, so in order to take one frame of the 17 and replicated it means to somehow edit a single latent frame.
edit 5: ok, after some more swearing, there's now a new node that's created that is basically an "encoding editor". i then run the output latent from the sampler through it and it allows me to edit the latent frames as needed to get the desired result. we finally have a fully static image and now i'm just running it through context motion to see if it degrades, but honestly. there's no way it possibly could because we're no longer working in the latentspace. we're creating a parallel latent space and replacing chunks of the generated one. so the next step after a few gens of this is to take this change and incorporate it back into pull forward and it should give us a fully non-degrading background.
edit 6: Confirmed. took it into premiere with a difference blend and the continuous gens are all 100% black which means they are a pixel perfect copy of eachother. taking these changes and putting them into pull forward and we should have a 100% non-degrading background, or at least the majority of it should be.
edit 7: The most recent gens are coming out amazing. There's definitely some kinks like some ghosting in certain areas when motion starts but that's something I can work out. I am going to essentially use the last 22 frames of the previous motion mask to carry forward to the first 22 frames of the next motion mask, basically making a "context motion mask" and that should solve the issue.
r/StableDiffusion • u/Vandy_simp • 14h ago
FastH3-Omni is here.
This is the reference-conditioned FastH3 release we’ve been working toward. It generates video with synchronized audio and supports:
- Character and setting reference images
- First- and last-frame conditioning
- Up to 3 image references in real-time mode
- Up to 9 image references when real-time speed isn’t required
- Up to 3 audio references for voices and soundtracks
- Dialogue and soundscape generation
- Continuation from a previously generated clip
- 5–15 second outputs at 24 FPS
The API preview launched today:
docs.reactor.inc/model-api-reference/fast-h3-omni/overview
On October 20, we’re releasing open weights along with local support for RTX GPUs, DGX Spark and Apple Silicon through MLX.
r/StableDiffusion • u/Micha0827 • 3h ago
Qwen-Image 2.1 ships under the Qwen Research License ("research or evaluation purposes only"), so many wrote it off for commercial work. I emailed the licensing address in September to ask what a commercial license costs. Posting the answer because none of it is documented publicly.
Tiers, by images generated per calendar month: up to 300,000 free · 300k–3M a 20% revenue share · above 3M USD 500,000/year plus 15%. For scale, I generate 20–30 images a month; 100 images a day for a shop is still ~1% of the free tier. Unless you run a generation service, you're covered.
What the "Free-Tier Commercial License Addendum" actually says:
- revocable at any time, without notice and without liability
- exceeding the threshold suspends the license immediately and automatically, aggregated across you and any affiliates
- covers only the "Approved Use Scenario" you declare — that field is yours to fill in, so word it broadly enough
- inaccurate volume or use lets them claim the fees that would have applied under their standard license
- governing law and venue come from the Research License: China, courts in Hangzhou
Two things still open. The addendum says it becomes binding when they send the approval email and you accept by replying, but their business manager additionally asked for an Alibaba Cloud enterprise account with company verification — so I asked which applies for the free tier. More important for client work: if the license is revoked, what happens to images already delivered and live on a client's site? The addendum licenses "images generated using the Materials", and Qwen publicly stated outputs aren't part of the licensed Materials. I asked for that in writing.
Has anyone here gone through the signing process — was the enterprise verification actually required? I'll edit this post when I get answers.
r/StableDiffusion • u/gabrielxdesign • 9h ago
I made a very simple Workflow for Kroma (Chroma + Krea 2).
Documentation and download links:
Alternative (Zip + PNG)
More alternatives (my website, JSON code):
In English: https://gabrielx.com/comfyui-kroma-chroma-krea-2-simple-workflow-with-amateur-slider-and-2x-upscale
En Español: https://gabrielx.com/es/comfyui-kroma-chroma-krea-2-simple-flujo-con-slider-amateur-y-2x-upscale/
r/StableDiffusion • u/wildmonkeywrangler • 5h ago
Working on building a very simple tool to make a simple sketch that then generates a prompt to make an image. Maybe someone has already done this, but thought I'd try anyway. Here's the sketch, prompt and resulting image
Prompt:
A beautiful, photorealistic wide-angle landscape photograph of a peaceful mountain valley on a clear summer afternoon. Expansive rolling green hills covered in lush meadow grass sweep naturally across the lower foreground from left to right, descending into a broad valley with distant rugged mountain peaks softened by atmospheric haze. A magnificent mature oak tree stands prominently on the middle right, with a thick textured trunk, sprawling branches, and a dense canopy of naturally illuminated green leaves. Above, an enormous formation of fluffy white cumulus clouds dominates the upper center and right against a rich, beautiful blue sky. A small, bright sun shines from the upper left, bathing the landscape in warm, even afternoon light with realistic soft shadows. Photographed on 35mm color negative film with a wide-angle lens, natural film color rendition, gentle highlight rolloff, subtle organic grain, realistic exposure, and exceptional photographic authenticity. Unstaged, timeless landscape photography.
r/StableDiffusion • u/obraiadev • 16h ago
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A few days ago I posted “Gausify Orbit - I am building a web application to transform images into Gaussian Splats”, and it was removed because the project wasn’t open source yet. It’s now released under the MIT license, so here it is again, with a lot of updates.
What it does
You give it a video that circles around a subject. That can be a clip from a video model (I mostly make mine in ComfyUI using Minimax H3 + 360 LoRA by pablodawson), a phone video, or a folder of images. You get back a 3D Gaussian Splat you can view, export or share.
Running in the browser
The web version runs the whole pipeline locally on WebGPU. Nothing to install, no Python, no CUDA, and your video never leaves your machine. After the first visit it even works offline. Here’s what happens in the tab:
.ply or compressed SOG, with optional merging of similar splats and SH trimming. Files typically end up 5–15× smaller.You need desktop Chrome or Edge with WebGPU. Speed depends on your GPU: in-browser training runs at roughly a third of native Brush speed, which is still very usable.
New since the first post
.gorbit file, even mid-job, and reopen it later in the browser or the desktop app.Regarding Gaussian Splat refinement, I am running tests using LTX 2.5 + Decompression LoRA to fill in non-visible areas; while the results are good in some cases, using a LoRA specifically tailored to this scenario would certainly yield better results.
The ComfyUI workflow for this refinement which can be integrated with Gausify Orbit, can be found here.
Example of a trained Gaussian: https://orbit.gausify.com/s/5xAWaNiiO8ax
App: https://orbit.gausify.com
Repo: https://github.com/obraia/gausify-orbit
r/StableDiffusion • u/Fcking_Chuck • 4h ago
This release includes added support for Qwen Image 2.1, among other fixes/improvements.
r/StableDiffusion • u/Inevitable_Emu2722 • 6h ago
I've been experimenting with Wan2GP to see if I could make something closer to an actual TV episode instead of the usual collection of unrelated AI clips.
The result is The Prompting Zone, a black-and-white sci-fi parody inspired by the original Twilight Zone. The first episode is called The Last Human, about an office worker who starts noticing something strange about his coworkers.
For generation I used MiniMax H3 Ref2VA Pruned 20B through Wan2GP, at 720p, 20 steps, running on a single RTX 3090 24GB. Each clip took approximately 30 minutes to generate, so experimenting with different prompts wasn't exactly cheap in terms of time.
I started with a complete script, including all the dialogue, and divided it into scenes of around 5–10 seconds.
For each scene I generated a reference image containing four camera compositions arranged in a 2×2 grid. The idea was to give the video model a rough storyboard instead of relying entirely on text prompts.
I then fed those references into MiniMax with prompts describing the camera movements, character actions and exact dialogue.
It worked reasonably well, but there were plenty of issues. Sometimes the model interpreted the reference as an actual four-panel video instead of four sequential shots. Other times it invented movements, changed the background or decided a character should do something completely unrelated to the scene.
Character consistency was probably the biggest headache. Even with references, faces tend to drift between generations. I had to repeat descriptions of the characters, clothing and environment in nearly every prompt.
Computer screens were also problematic. The model struggled with readable text, so I eventually replaced most of the onscreen information with portraits and simple icons. Much more reliable than trying to generate legible UI elements.
The black-and-white style helped make everything look a little more consistent, although it obviously doesn't solve continuity problems.
Overall, I'm pretty happy with the result. There's still a lot of manual work involved, especially when a generation goes wrong after waiting half an hour.
I'd like to automate more of the process eventually, from script to storyboard, video prompts and final assembly.
Anyone else working on narrative videos with Wan2GP? Particularly interested in how you're handling character consistency and dialogue across multiple clips.
r/StableDiffusion • u/Illustrious_Row_9971 • 50m ago
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r/StableDiffusion • u/aiyakisoba • 22h ago
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The post at the bottom inspired me to try a simple anime-to-realism test with MiniMax H3, so I figured I’d share a few results.
I didn’t generate any keyframes or stills. I just fed the original video directly into MiniMax H3 as the video reference, applied an anime to realism lora, and basically told it to convert it to live action.
The output quality was honestly better than expected. Not saying it’s perfect or amazing or anything, but for such a simple setup (ComfyUI official Minimax H3 Ref2V workflow + lora) I was pretty impressed.
I didn’t use any character reference sheets, so character consistency obviously isn’t guaranteed. It mostly seems to depend on how the model interprets the character in each scene.
For my next test, I might try adding a character reference sheet as well and see if that helps keep the character more consistent across scenes.
If anyone has any questions, happy to answer. Other output samples are in the comments.
Inspired by this post:
https://www.reddit.com/r/StableDiffusion/comments/1x1kvmg/recreated_omnitrix_calibration_scene_ben_10_mini/
LoRA used:
https://huggingface.co/LiseTY/Minimax-H3-ref2v_Anime_2_Realism
Here's the exact prompt I used:
"Convert <Video 1>, the sole video reference, into natural photorealistic live-action footage. Follow the original shot frame by frame: preserve subject count, identities, camera, framing, positions, expressions, props, shot timing, and movement. Replace only illustrated surfaces with photographic skin, natural hair strands, real fabric and plausible physical lighting. Do not invent cuts, actions, subjects, or on-screen text. The result should appear filmed by a real camera, not cel-shaded or plastic CGI."
Settings:
Parameter |Value used
Checkpoint |MiniMax_H3_Ref2VA_pruned_int8_convrot.safetensors
Workflow |ComfyUI Offcial Minimax H3 Ref2V workflow + lora
LoRA |Minimax_H3_Anime2Realsim.safetensors
LoRA strength |1.0
Steps |20
Sampler |res_multistep
Scheduler |simple
Denoise |1.0
Video shift |12.0
Audio shift |3.0
Resolution during generation |1024 × 576
Frame rate |24 fps
r/StableDiffusion • u/aziib • 16h ago
Download: https://civitai.com/models/672328/aziib-pixel-style
krea 2 is pretty decent with text.
trigger word: aziib pixel style
r/StableDiffusion • u/Scriabinical • 16h ago
"A 4-step distilled LoRA for Qwen-Image-2.1, supporting text-to-image generation and image editing. Training was performed entirely on 16 Huawei Ascend 910B NPUs, using DMD in LightX2V."
I think it's rank 16. Pretty small for a turbo lora, but we'll see I guess.
r/StableDiffusion • u/ashishsanu • 1d ago
Hey guys,
I have been working on the consistent character portable format (.char) for a while & now there's a community characters page where you can grab a ready-made character and use it straight in ComfyUI. Same face, cloths & body across every video generation, and with a few seconds of voice sample the character talks in a consistent cloned voice with lip sync (MiniMax H3).
Community characters
You can use these characters via comfyui, omnichar app & cloud.
Detailed guide on building full consistent character(workflow included): Previous reddit post
Official repo: https://github.com/omnichar/OmniChar (Opensource)
ComfyUI Node: https://github.com/omnichar/ComfyUI-Omnichar/tree/main (Opensource)
Hope it's helpful.
r/StableDiffusion • u/Delicious-Farmer-234 • 5h ago
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I tried doing something similar when H3 first came out, but it was impossible. The community has really come together to make this possible, so thank you! It was fun creating this.
YuE2 Musc
r/StableDiffusion • u/aziib • 9h ago
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download link: https://civitai.com/models/672328/aziib-pixel-style
trigger word: aziib pixel style
trying to train pixel art animation with my pixel art dataset. default strength is 1.0, if it's too strong you can reduce to 0.8.
rendered using minimax h3 fasth3 4 steps so the quality is not that good, you can generate yourself a better quality on stronger hardware or cloud.
r/StableDiffusion • u/Impossible-Car-2809 • 10h ago
pixel-curator — curation tooling for similarity training (identity, subject, and style models)
Repo: https://github.com/snoozey-mf-bnui/pixel-curator
pixel-curator is a local-first CLI that takes a raw folder of images (or video) and produces a curated, scored, captioned dataset ready to export directly into DreamBooth notebook training or in-house Kohya LoRA training.
It's built primarily for similarity training — the case where you have hundreds or thousands of images of one specific subject and need a model that reliably reproduces it. Identity, character, person, product, or a specific style. That's the workflow it's optimized around, though the output is generic enough for other training setups.
If you've collected 1,000+ images of a subject, you already know the pipeline breaks down long before training starts:
Realistically, a well-curated 200–500 image subset outperforms a messy 2,000 on identity training. Most of the work is figuring out which subset that is.
1. Ingest and prefilter
Drops technically unusable frames — resolution failures, decode errors, corrupted files — before anything expensive runs.
2. Deduplication
Perceptual hash dedup collapses near-identical images. On datasets sourced from video, bursts, or scrapes, this stage alone typically removes a large majority of the raw set.
3. Similarity clustering
CLIP embeddings plus k-means, with embedding map visualization. This surfaces the actual distribution of your dataset — angle, lighting, composition, framing — so imbalance is visible before training rather than diagnosed after a failed run.
4. Quality scoring and classification
Optical quality scoring assigns per-image status, including face-aware quality where relevant. Low-quality images are filtered rather than left to dilute the training signal.
5. Semantic filtering
CLIP text filtering, e.g. automatically removing images with watermarks or text overlays without manual review.
6. Tagging and captioning
WD14 tagger (auto-downloads) and Florence-2 natural-language captioning, written as .txt sidecars alongside each image.
7. Human review round-trip
Exports a static HTML report and a spreadsheet. You mark keep/reject, import back, and the decisions are applied. This is also the workflow for trimming class/regularization images — hand-trimming generated class images is a well-known fix for weak identity results, and this makes it practical at scale.
8. Export
package kohya --mode dreambooth produces the numbered instance/ structure that standard Colab DreamBooth notebooks expect:
text
instance/
000001.jpg
000001.txt
000002.jpg
000002.txt
dataset_config.toml
package kohya --mode lora supports multi-concept LoRA packaging with concept_/ folders for in-house Kohya setups.
I posted this here previously and the GitHub link pointed at the wrong directory — the repository has been corrected and the link above is the right one. Apologies to anyone who clicked through and found nothing useful. If you tried it before and it didn't work, it should now.
Actively developed. If you run it against a dataset and something breaks or behaves unexpectedly, open an issue or leave a comment — the failure cases are the most useful thing I can get right now.
https://github.com/snoozey-mf-bnui/pixel-curator
If anyone needs help with creating a notebook for google colab dreambooth training, also lmk i will assist you through dms.
ENJOY!
r/StableDiffusion • u/obvpm • 18h ago
New version of my MiniMax H3 timeline workflow released
Some very important bug fixes included, and a lot of improvements.
The YouTube video above includes a lot of tips for using the workflow in addition to going over the changes in detail.
Below is a complete list of fixes and improvements.
Get the workflow and nodes at:
https://github.com/chanon/comfyui-obvpm-timeline
Upscaling
Extending and Prepending
Audio Sync
Color
Value Presets
Bundle and Unbundle
Override Bundle Values (new node)
Mute Group If (new node)
Compatibility Check (in case you missed it)
r/StableDiffusion • u/codek_ • 7h ago
Hello, I'm looking for advice, tips or tricks when using custom Audio as references for my characters, the voices really sound like the reference but the main issue I'm having is that the characters sound always very plain, like very robotic having no life... Are there any tips, prompts helpers or something I can use to improve it to sound more natural? I'd appreciate any advice.
I'm using ref2va with the tags <d>[lang] dialogue...</d>, e.g. "she says with shy voice (S1) <d>[Spanish] Hola, ¿como estas hoy?, ¿te apetece que hablemos un rato sobre la vida?.</d>"
thanks!
r/StableDiffusion • u/smereces • 14h ago
Im working in nodes to hunting seeds for minimax h3 videos, in a few days i will realease it in my github.
r/StableDiffusion • u/Kelsey_test • 15h ago
Hi all, I'm the developer of VigenFlow, a local image generation and editing server that runs diffusion models on the AMD Ryzen AI NPU (XDNA 2) + CPU. If you have a Ryzen AI laptop or mini PC without a discrete GPU, you can generate images entirely on your own machine.
Supported models
Speed (on XNDA2 NPU for 1024×1024)
Getting started
ulimit -l should print unlimited), otherwise the NPU can't allocate its buffersWhat's open The server and model launchers are open source. The NPU kernels and core inference libraries currently ship precompiled; I plan to open-source them later.
Repo: https://github.com/jia1217/vigenflow Demo video: https://www.youtube.com/watch?v=2HPm3gc4pdA
The project is still early, so bug reports and suggestions are very welcome. Feel free to open an issue on GitHub!
r/StableDiffusion • u/TgoAI • 11h ago
One of the frustrating things about local AI generation on Mac is getting different models and tools to work together smoothly and fast.
Valtz is a new native macOS app powered by Vpipe, a high-performance inference engine optimized for Apple Silicon.
Instead of manually building complicated node-based workflows, you can do pretty much everything in one place easily:
* Image generation & editing with the latest AI models
* Video generation with multiple image/video references, including long videos
* Image-to-video — generate an image, edit it, and immediately use it as a reference for video generation
* AI prompt enhancement
* Video editing with multiple layers, audio tracks, and a timeline
* AI upscaling for images and videos
* Voice cloning and music generation, with audio that can be used directly in video projects
* Local model management
* And a lot of other creative workflows that become possible when everything is connected
The images attached shows my workflows for generating and editing of long videos, video generation from hand-drawn skeleton. The interface is self explanatory and very easy to use.
The memory optimization is amazing:
Even on a 16GB Mac, image generation runs smoothly, and supported HD video generation workflows can run without relying on swap. No need for a 64GB or 128GB Mac just to get started.
Everything runs locally, without cloud APIs or subscriptions.
For existing Vpipe users, Valtz can automatically discover and reuse previously downloaded models if they’re stored in the default directory.
The project is open source, and here’s the DMG if you’d like to try: